General fine tuning for Stable Diffusion
Go to file
Xavier 0c9611f251
Update README.md
2022-09-06 09:35:58 -07:00
configs change lr 2022-09-06 09:32:32 -07:00
evaluation initial commit 2022-09-06 00:00:21 -07:00
img initial commit 2022-09-06 00:00:21 -07:00
ldm initial commit 2022-09-06 00:00:21 -07:00
models initial commit 2022-09-06 00:00:21 -07:00
scripts initial commit 2022-09-06 00:00:21 -07:00
LICENSE initial commit 2022-09-06 00:00:21 -07:00
README.md Update README.md 2022-09-06 09:35:58 -07:00
environment.yaml initial commit 2022-09-06 00:00:21 -07:00
main.py initial commit 2022-09-06 00:00:21 -07:00
merge_embeddings.py initial commit 2022-09-06 00:00:21 -07:00
setup.py initial commit 2022-09-06 00:00:21 -07:00

README.md

Dreambooth on Stable Diffusion

This is an implementtaion of Google's Dreambooth with Stable Diffusion. The original Dreambooth is based on Imagen text-to-image model. However, neither the model nor the pre-trained weights of Imagen is available. To enable people to fine-tune a text-to-image model with a few examples, I implemented the idea of Dreambooth on Stable diffusion.

This code repository is based on that of Textual Inversion. Note that Textual Inversion only optimizes word ebedding, while dreambooth fine-tunes the whole diffusion model.

The implementation makes minimum changes over the official codebase of Textual Inversion. In fact, due to lazyness, some components in Textual Inversion, such as the embedding manager, are not deleted, although they will never be used here.

Usage

Preparation

To fine-tune a stable diffusion model, you need to obtain the pre-trained stable diffusion models following their instructions. Weights can be downloads on HuggingFace. You can decide which version of checkpoint to use, but I use sd-v1-4-full-ema.ckpt.

We also need to create a set of images for regularization, as the fine-tuning algorithm of Dreambooth requires that. Details of the algorithm can be found in the paper. The text prompt can be photo of a xxx, where xxx is a word that describes the class of your object, such as dog. The command is

python scripts/stable_txt2img.py --ddim_eta 0.0 --n_samples 8 --n_iter 1 --scale 10.0 --ddim_steps 50  --ckpt /path/to/original/stable-diffusion/sd-v1-4-full-ema.ckpt --prompt "a photo of a <xxx>" 

I generate 8 images for regularization. After that, save the generated images (separately, one image per .png file) at /root/to/regularization/images.

Training

Training can be done by running the following command

python main.py --base configs/stable-diffusion/v1-finetune_unfrozen.yaml 
                -t 
                --actual_resume /path/to/original/stable-diffusion/sd-v1-4-full-ema.ckpt  
                -n <job name> 
                --gpus 0, 
                --data_root /root/to/training/images 
                --reg_data_root /root/to/regularization/images 
                --class_word <xxx>

Detailed configuration can be found in configs/stable-diffusion/v1-finetune_unfrozen.yaml. In particular, the default learning rate is 1.0e-6 as I found the 1.0e-5 in the Dreambooth paper leads to poor editability. The parameter reg_weight corresponds to the weight of regularization in the Dreambooth paper, and the default is set to 1.0.

Generation