support for sdxl-inpaint model
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model:
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target: sgm.models.diffusion.DiffusionEngine
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params:
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scale_factor: 0.13025
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disable_first_stage_autocast: True
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denoiser_config:
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target: sgm.modules.diffusionmodules.denoiser.DiscreteDenoiser
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params:
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num_idx: 1000
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weighting_config:
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target: sgm.modules.diffusionmodules.denoiser_weighting.EpsWeighting
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scaling_config:
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target: sgm.modules.diffusionmodules.denoiser_scaling.EpsScaling
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discretization_config:
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target: sgm.modules.diffusionmodules.discretizer.LegacyDDPMDiscretization
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network_config:
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target: sgm.modules.diffusionmodules.openaimodel.UNetModel
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params:
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adm_in_channels: 2816
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num_classes: sequential
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use_checkpoint: True
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in_channels: 9
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out_channels: 4
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model_channels: 320
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attention_resolutions: [4, 2]
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num_res_blocks: 2
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channel_mult: [1, 2, 4]
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num_head_channels: 64
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use_spatial_transformer: True
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use_linear_in_transformer: True
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transformer_depth: [1, 2, 10] # note: the first is unused (due to attn_res starting at 2) 32, 16, 8 --> 64, 32, 16
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context_dim: 2048
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spatial_transformer_attn_type: softmax-xformers
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legacy: False
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conditioner_config:
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target: sgm.modules.GeneralConditioner
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params:
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emb_models:
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# crossattn cond
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- is_trainable: False
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input_key: txt
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target: sgm.modules.encoders.modules.FrozenCLIPEmbedder
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params:
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layer: hidden
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layer_idx: 11
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# crossattn and vector cond
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- is_trainable: False
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input_key: txt
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target: sgm.modules.encoders.modules.FrozenOpenCLIPEmbedder2
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params:
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arch: ViT-bigG-14
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version: laion2b_s39b_b160k
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freeze: True
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layer: penultimate
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always_return_pooled: True
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legacy: False
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# vector cond
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- is_trainable: False
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input_key: original_size_as_tuple
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target: sgm.modules.encoders.modules.ConcatTimestepEmbedderND
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params:
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outdim: 256 # multiplied by two
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# vector cond
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- is_trainable: False
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input_key: crop_coords_top_left
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target: sgm.modules.encoders.modules.ConcatTimestepEmbedderND
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params:
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outdim: 256 # multiplied by two
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# vector cond
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- is_trainable: False
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input_key: target_size_as_tuple
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target: sgm.modules.encoders.modules.ConcatTimestepEmbedderND
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params:
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outdim: 256 # multiplied by two
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first_stage_config:
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target: sgm.models.autoencoder.AutoencoderKLInferenceWrapper
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params:
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embed_dim: 4
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monitor: val/rec_loss
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ddconfig:
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attn_type: vanilla-xformers
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double_z: true
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z_channels: 4
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resolution: 256
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in_channels: 3
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out_ch: 3
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ch: 128
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ch_mult: [1, 2, 4, 4]
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num_res_blocks: 2
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attn_resolutions: []
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dropout: 0.0
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lossconfig:
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target: torch.nn.Identity
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@ -106,6 +106,20 @@ def txt2img_image_conditioning(sd_model, x, width, height):
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return x.new_zeros(x.shape[0], 2*sd_model.noise_augmentor.time_embed.dim, dtype=x.dtype, device=x.device)
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else:
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sd = sd_model.model.state_dict()
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diffusion_model_input = sd.get('diffusion_model.input_blocks.0.0.weight', None)
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if diffusion_model_input.shape[1] == 9:
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# The "masked-image" in this case will just be all 0.5 since the entire image is masked.
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image_conditioning = torch.ones(x.shape[0], 3, height, width, device=x.device) * 0.5
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image_conditioning = images_tensor_to_samples(image_conditioning,
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approximation_indexes.get(opts.sd_vae_encode_method))
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# Add the fake full 1s mask to the first dimension.
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image_conditioning = torch.nn.functional.pad(image_conditioning, (0, 0, 0, 0, 1, 0), value=1.0)
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image_conditioning = image_conditioning.to(x.dtype)
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return image_conditioning
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# Dummy zero conditioning if we're not using inpainting or unclip models.
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# Still takes up a bit of memory, but no encoder call.
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# Pretty sure we can just make this a 1x1 image since its not going to be used besides its batch size.
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@ -362,6 +376,11 @@ class StableDiffusionProcessing:
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if self.sampler.conditioning_key == "crossattn-adm":
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return self.unclip_image_conditioning(source_image)
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sd = self.sampler.model_wrap.inner_model.model.state_dict()
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diffusion_model_input = sd.get('diffusion_model.input_blocks.0.0.weight', None)
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if diffusion_model_input.shape[1] == 9:
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return self.inpainting_image_conditioning(source_image, latent_image, image_mask=image_mask)
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# Dummy zero conditioning if we're not using inpainting or depth model.
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return latent_image.new_zeros(latent_image.shape[0], 5, 1, 1)
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@ -15,6 +15,7 @@ config_sd2v = os.path.join(sd_repo_configs_path, "v2-inference-v.yaml")
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config_sd2_inpainting = os.path.join(sd_repo_configs_path, "v2-inpainting-inference.yaml")
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config_sdxl = os.path.join(sd_xl_repo_configs_path, "sd_xl_base.yaml")
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config_sdxl_refiner = os.path.join(sd_xl_repo_configs_path, "sd_xl_refiner.yaml")
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config_sdxl_inpainting = os.path.join(sd_configs_path, "sd_xl_inpaint.yaml")
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config_depth_model = os.path.join(sd_repo_configs_path, "v2-midas-inference.yaml")
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config_unclip = os.path.join(sd_repo_configs_path, "v2-1-stable-unclip-l-inference.yaml")
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config_unopenclip = os.path.join(sd_repo_configs_path, "v2-1-stable-unclip-h-inference.yaml")
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@ -71,7 +72,10 @@ def guess_model_config_from_state_dict(sd, filename):
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sd2_variations_weight = sd.get('embedder.model.ln_final.weight', None)
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if sd.get('conditioner.embedders.1.model.ln_final.weight', None) is not None:
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return config_sdxl
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if diffusion_model_input.shape[1] == 9:
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return config_sdxl_inpainting
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else:
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return config_sdxl
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if sd.get('conditioner.embedders.0.model.ln_final.weight', None) is not None:
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return config_sdxl_refiner
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elif sd.get('depth_model.model.pretrained.act_postprocess3.0.project.0.bias', None) is not None:
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@ -34,6 +34,11 @@ def get_learned_conditioning(self: sgm.models.diffusion.DiffusionEngine, batch:
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def apply_model(self: sgm.models.diffusion.DiffusionEngine, x, t, cond):
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sd = self.model.state_dict()
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diffusion_model_input = sd.get('diffusion_model.input_blocks.0.0.weight', None)
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if diffusion_model_input.shape[1] == 9:
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x = torch.cat([x] + cond['c_concat'], dim=1)
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return self.model(x, t, cond)
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