From e1367298c7497de7b9a5f076f40433ce6857077e Mon Sep 17 00:00:00 2001 From: apolinario Date: Tue, 16 Aug 2022 21:08:28 +0200 Subject: [PATCH] Add diffusers as a way to inference in the model --- README.md | 19 +++++++++++++++++++ 1 file changed, 19 insertions(+) diff --git a/README.md b/README.md index 63e96be..4390916 100644 --- a/README.md +++ b/README.md @@ -87,6 +87,25 @@ and sample with ``` python scripts/txt2img.py --prompt "a photograph of an astronaut riding a horse" --plms ``` + +Another way to download and sample Stable Diffusion is by using the [diffusers library](https://github.com/huggingface/diffusers/tree/main#new--stable-diffusion-is-now-fully-compatible-with-diffusers) +```py +# make sure you're logged in with `huggingface-cli login` +from torch import autocast +from diffusers import StableDiffusionPipeline, LMSDiscreteScheduler + +pipe = StableDiffusionPipeline.from_pretrained( + "CompVis/stable-diffusion-v1-3-diffusers", + use_auth_token=True +) + +prompt = "a photo of an astronaut riding a horse on mars" +with autocast("cuda"): + image = pipe(prompt)["sample"][0] + +image.save("astronaut_rides_horse.png") +``` + By default, this uses a guidance scale of `--scale 7.5`, [Katherine Crowson's implementation](https://github.com/CompVis/latent-diffusion/pull/51) of the [PLMS](https://arxiv.org/abs/2202.09778) sampler, and renders images of size 512x512 (which it was trained on) in 50 steps. All supported arguments are listed below (type `python scripts/txt2img.py --help`).