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README.md
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README.md
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@ -78,6 +78,9 @@ steps show the relative improvements of the checkpoints:
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Stable Diffusion is a latent diffusion model conditioned on the (non-pooled) text embeddings of a CLIP ViT-L/14 text encoder.
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#### Sampling Script
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After [obtaining the weights](#weights), link them
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```
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mkdir -p models/ldm/stable-diffusion-v1/
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@ -88,24 +91,6 @@ and sample with
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python scripts/txt2img.py --prompt "a photograph of an astronaut riding a horse" --plms
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```
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Another way to download and sample Stable Diffusion is by using the [diffusers library](https://github.com/huggingface/diffusers/tree/main#new--stable-diffusion-is-now-fully-compatible-with-diffusers)
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```py
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# make sure you're logged in with `huggingface-cli login`
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from torch import autocast
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from diffusers import StableDiffusionPipeline, LMSDiscreteScheduler
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pipe = StableDiffusionPipeline.from_pretrained(
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"CompVis/stable-diffusion-v1-3-diffusers",
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use_auth_token=True
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)
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prompt = "a photo of an astronaut riding a horse on mars"
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with autocast("cuda"):
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image = pipe(prompt)["sample"][0]
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image.save("astronaut_rides_horse.png")
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```
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By default, this uses a guidance scale of `--scale 7.5`, [Katherine Crowson's implementation](https://github.com/CompVis/latent-diffusion/pull/51) of the [PLMS](https://arxiv.org/abs/2202.09778) sampler,
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and renders images of size 512x512 (which it was trained on) in 50 steps. All supported arguments are listed below (type `python scripts/txt2img.py --help`).
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@ -149,6 +134,28 @@ non-EMA to EMA weights. If you want to examine the effect of EMA vs no EMA, we p
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which contain both types of weights. For these, `use_ema=False` will load and use the non-EMA weights.
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#### Diffusers Integration
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Another way to download and sample Stable Diffusion is by using the [diffusers library](https://github.com/huggingface/diffusers/tree/main#new--stable-diffusion-is-now-fully-compatible-with-diffusers)
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```py
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# make sure you're logged in with `huggingface-cli login`
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from torch import autocast
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from diffusers import StableDiffusionPipeline, LMSDiscreteScheduler
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pipe = StableDiffusionPipeline.from_pretrained(
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"CompVis/stable-diffusion-v1-3-diffusers",
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use_auth_token=True
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)
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prompt = "a photo of an astronaut riding a horse on mars"
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with autocast("cuda"):
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image = pipe(prompt)["sample"][0]
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image.save("astronaut_rides_horse.png")
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```
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### Image Modification with Stable Diffusion
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By using a diffusion-denoising mechanism as first proposed by [SDEdit](https://arxiv.org/abs/2108.01073), the model can be used for different
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